stable diffusion v1 5

Providerstable-diffusion-v1-5
Categorytext-to-image
Licensecreativeml-openrail-m
Downloads1.5M
Stars0

Overview

Stable Diffusion v1.5 remains a foundational pillar for open-source generative AI, offering a versatile text-to-image pipeline that balances performance with hardware accessibility. Unlike closed-API models, v1.5 is designed for local deployment and deep customization, making it the primary target for community-driven fine-tuning. Developers can leverage its latent diffusion architecture to implement custom LoRAs or ControlNets, granting precise structural control over image generation that exceeds basic prompting. Whether you are building an automated asset pipeline, an AI-powered design tool, or integrating image synthesis into a full-stack app, v1.5 provides a stable, well-documented baseline with massive ecosystem support and low VRAM overhead compared to newer, larger models.

Highlights

  • Open-weights architecture allowing for local deployment and hosting
  • Extensive ecosystem support for LoRA and ControlNet extensions
  • Low VRAM requirements suitable for consumer-grade GPU hardware
  • Highly flexible integration via Diffusers and ComfyUI frameworks
  • Proven reliability for rapid prototyping of image-gen pipelines

Usage

Install
# Install Hugging Face transformers
pip install transformers torch
SDK Usage
# Load model with transformers
from transformers import AutoModel, AutoTokenizer

model = AutoModel.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5")
tokenizer = AutoTokenizer.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5")

Hugging Face Download

We recommend downloading the model via the Hugging Face CLI or Hub SDK.

Guidance:Before downloading, install huggingface_hub with:

Guidance
pip install -U huggingface_hub

CLI Download

Download the full repository

Download the full repository
huggingface-cli download stable-diffusion-v1-5/stable-diffusion-v1-5

Download a single file to a local folder (e.g. config.json into ./dir)

Download a single file to a local folder (e.g. config.json into ./dir)
huggingface-cli download stable-diffusion-v1-5/stable-diffusion-v1-5 config.json --local-dir ./dir

See the official docs for more CLI options

SDK Download

SDK Download
# 模型下载
from huggingface_hub import snapshot_download
model_dir = snapshot_download('stable-diffusion-v1-5/stable-diffusion-v1-5')

Git Download

Make sure git-lfs is installed first

Git Download
git lfs install
git clone https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5

To skip LFS large-file downloads, use:

Skip LFS
GIT_LFS_SKIP_SMUDGE=1 git clone https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5

Model files are hosted on the Hugging Face Hub — download directly via HF CLI / SDK / Git, not through this site.

PyTorch / Transformers Usage

Install Transformers

Install Transformers
pip install -U transformers torch

Load the model and run inference

Load the model and run inference
from transformers import AutoModelForCausalLM, AutoTokenizer

model = AutoModelForCausalLM.from_pretrained('stable-diffusion-v1-5/stable-diffusion-v1-5')
tokenizer = AutoTokenizer.from_pretrained('stable-diffusion-v1-5/stable-diffusion-v1-5')

Full Documentation

来源: HuggingFace

---
license: creativeml-openrail-m
tags:

  • stable-diffusion

  • stable-diffusion-diffusers

  • text-to-image

inference: true
---

Stable Diffusion v1-5 Model Card

⚠️ This repository is a mirror of the now deprecated ruwnayml/stable-diffusion-v1-5, this repository or organization are not affiliated in any way with RunwayML.

Modifications to the original model card are in <span style="color:crimson">red</span> or <span style="color:darkgreen">green</span>

Stable Diffusion is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input.
For more information about how Stable Diffusion functions, please have a look at 🤗's Stable Diffusion blog.

The Stable-Diffusion-v1-5 checkpoint was initialized with the weights of the Stable-Diffusion-v1-2
checkpoint and subsequently fine-tuned on 595k steps at resolution 512x512 on "laion-aesthetics v2 5+" and 10% dropping of the text-conditioning to improve classifier-free guidance sampling.

You can use this both with the 🧨Diffusers library and RunwayML GitHub repository (<span style="color:crimson">now deprecated</span>), <span style="color:darkgreen">ComfyUI, Automatic1111, SD.Next, InvokeAI</span>.

Use with Diffusers

py
from diffusers import StableDiffusionPipeline
import torch

model_id = "sd-legacy/stable-diffusion-v1-5"
pipe = StableDiffusionPipeline.from_pretrained(model_id, torch_dtype=torch.float16)
pipe = pipe.to("cuda")

prompt = "a photo of an astronaut riding a horse on mars"
image = pipe(prompt).images[0]

image.save("astronaut_rides_horse.png")


For more detailed instructions, use-cases and examples in JAX follow the instructions here

Use with GitHub Repository <span style="color:crimson">(now deprecated)</span>, <span style="color:darkgreen">ComfyUI or Automatic1111</span>

1. Download the weights
- v1-5-pruned-emaonly.safetensors - ema-only weight. uses less VRAM - suitable for inference
- v1-5-pruned.safetensors - ema+non-ema weights. uses more VRAM - suitable for fine-tuning

2. Follow instructions here. <span style="color:crimson">(now deprecated)</span>

3. <span style="color:darkgreen">Use locally with <a href="https://github.com/comfyanonymous/ComfyUI">ComfyUI</a>, <a href="https://github.com/AUTOMATIC1111/stable-diffusion-webui">AUTOMATIC1111</a>, <a href="https://github.com/vladmandic/automatic">SD.Next</a>, <a href="https://github.com/invoke-ai/InvokeAI">InvokeAI</a></span>

Model Details

  • Developed by: Robin Rombach, Patrick Esser
  • Model type: Diffusion-based text-to-image generation model
  • Language(s): English
  • Cite as:

@InProceedings{Rombach_2022_CVPR,
author = {Rombach, Robin and Blattmann, Andreas and Lorenz, Dominik and Esser, Patrick and Ommer, Bj\"orn},
title = {High-Resolution Image Synthesis With Latent Diffusion Models},
booktitle = {Proceedings of the IEEE/CVF Conference on Computer Vision and Pattern Recognition (CVPR)},
month = {June},
year = {2022},
pages = {10684-10695}
}

Uses

Direct Use

The model is intended for research purposes only. Possible research areas and tasks include
  • Safe deployment of models which have the potential to generate harmful content.
  • Probing and understanding the limitations and biases of generative models.
  • Generation of artworks and use in design and other artistic processes.
  • Applications in educational or creative tools.
  • Research on generative models.

Excluded uses are described below.

### Misuse, Malicious Use, and Out-of-Scope Use
_Note: This section is taken from the DALLE-MINI model card, but applies in the same way to Stable Diffusion v1_.

The model should not be used to intentionally create or disseminate images that create hostile or alienating environments for people. This includes generating images that people would foreseeably find disturbing, distressing, or offensive; or content that propagates historical or current stereotypes.

#### Out-of-Scope Use
The model was not trained to be factual or true representations of people or events, and therefore using the model to generate such content is out-of-scope for the abilities of this model.

#### Misuse and Malicious Use
Using the model to generate content that is cruel to individuals is a misuse of this model. This includes, but is not limited to:

  • Generating demeaning, dehumanizing, or otherwise harmful representations of people or their environments, cultures, religions, etc.
  • Intentionally promoting or propagating discriminatory content or harmful stereotypes.
  • Impersonating individuals without their consent.
  • Sexual content without consent of the people who might see it.
  • Mis- and disinformation
  • Representations of egregious violence and gore
  • Sharing of copyrighted or licensed material in violation of its terms of use.
  • Sharing content that is an alteration of copyrighted or licensed material in violation of its terms of use.

Limitations and Bias

Limitations

  • The model does not achieve perfect photorealism
  • The model cannot render legible text
  • The model does not perform well on more difficult tasks which involve compositionality, such as rendering an image corresponding to “A red cube on top of a blue sphere”
  • Faces and people in general may not be generated properly.
  • The model was trained mainly with English captions and will not work as well in other languages.
  • The autoencoding part of the model is lossy
  • The model was trained on a large-scale dataset
LAION-5B which contains adult material and is not fit for product use without additional safety mechanisms and considerations.
  • No additional measures were used to deduplicate the dataset. As a result, we observe some degree of memorization for images that are duplicated in the training data.
The training data can be searched at https://rom1504.github.io/clip-retrieval/ to possibly assist in the detection of memorized images.

Bias

While the capabilities of image generation models are impressive, they can also reinforce or exacerbate social biases.
Stable Diffusion v1 was trained on subsets of [LAION-2B(en)](https://laion.ai/blog/laion-5b

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